Edit an existing image to fit a given text description

powered by Stable Diffusion AI (CreativeML Open RAIL-M)

Prompt

Describe how the final image should look like

Strength

How strongly the original image should be altered (from subtle to drastic changes)

80%

Init image

Upload the original image to be modified.

Max: 512x512px (auto-resized)

Drag and drop files here or click

Model

The AI used to generate the image.

DreamShaper 8 (general)

Available in Power Mode

Count

Number of images to generate

More options

LoRA

Choose extension models

Setup LoRA (0 enabled)

Negative prompt

Describe what you DON'T want in the generated image

ugly, tiling, poorly drawn hands, poorly drawn feet, poorly drawn face, out of frame, extra limbs, disfigured, deformed, body out of frame, blurry, bad anatomy, blurred, watermark, grainy, signature, cut off, draft

Guidance

Adjusts how much the AI tries to fit the prompt (higher = stricter, lower = more freedom). The sweet spot is between 6-10, extreme values may produce more artifacts.

7

Available in Power Mode

Steps

Number of sampling steps. More steps = more details but also longer computation time.

Sampler

Defines the sampling method used to generate the image.

DPM++ 2M Karras

Seed

Unique image seed number. If not provided, the image will be random.

Available in Power Mode

Upscale

Upscale the image by this factor using the Real-ESRGAN model. Currently only a factor of 2 is supported.

Format

The encoding format of the generated image.

PNG